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Fine-tune FLUX fast
Customize FLUX.1 [dev] with the fast FLUX trainer on Replicate
Train the model to recognize and generate new concepts using a small set of example images, for specific styles, characters, or objects. It's fast (under 2 minutes), cheap (under $2), and gives you a warm, runnable model plus LoRA weights to download.
Featured models
google / veo-3
Sound on: Google’s flagship Veo 3 text to video model, with audio
luma / reframe-video
Change the aspect ratio of any video up to 30 seconds long, outputs will be 720p
google / imagen-4-fast
Use this fast version of Imagen 4 when speed and cost are more important than quality
google / imagen-4-ultra
Use this ultra version of Imagen 4 when quality matters more than speed and cost
google / imagen-4
Google's Imagen 4 flagship model
replicate / fast-flux-trainer
Train subjects or styles faster than ever
black-forest-labs / flux-kontext-pro
A state-of-the-art text-based image editing model that delivers high-quality outputs with excellent prompt following and consistent results for transforming images through natural language
black-forest-labs / flux-kontext-max
A premium text-based image editing model that delivers maximum performance and improved typography generation for transforming images through natural language prompts
ideogram-ai / ideogram-v3-turbo
Turbo is the fastest and cheapest Ideogram v3. v3 creates images with stunning realism, creative designs, and consistent styles
Official models
Official models are always on, maintained, and have predictable pricing.
I want to…
Generate images
Models that generate images from text prompts
Make videos with Wan2.1
Generate videos with Wan2.1, the fastest and highest quality open-source video generation model.
Generate videos
Models that create and edit videos
Caption images
Models that generate text from images
Transcribe speech
Models that convert speech to text
Use the FLUX family of models
The FLUX family of text-to-image models from Black Forest Labs
Remove backgrounds
Models that remove backgrounds from images and videos
Restore images
Models that improve or restore images by deblurring, colorization, and removing noise
Caption videos
Models that generate text from videos
Edit images
Tools for manipulating images.
Use a face to make images
Make realistic images of people instantly
Get embeddings
Models that generate embeddings from inputs
Generate speech
Convert text to speech
Generate music
Models to generate and modify music
Generate text
Models that can understand and generate text
Use handy tools
Toolbelt-type models for videos and images.
Upscale images
Upscaling models that create high-quality images from low-quality images
Use official models
Official models are always on, maintained, and have predictable pricing.
Enhance videos
Models that enhance videos with super-resolution, sound effects, motion capture and other useful production effects.
Detect objects
Models that detect or segment objects in images and videos.
Sing with voices
Voice-to-voice cloning and musical prosody
Make 3D stuff
Models that generate 3D objects, scenes, radiance fields, textures and multi-views.
Chat with images
Ask language models about images
Extract text from images
Optical character recognition (OCR) and text extraction
Use FLUX fine-tunes
Browse the diverse range of fine-tunes the community has custom-trained on Replicate
Control image generation
Guide image generation with more than just text. Use edge detection, depth maps, and sketches to get the results you want.
Popular models
SDXL-Lightning by ByteDance: a fast text-to-image model that makes high-quality images in 4 steps
This is an optimised version of the FLUX.1 [schnell] model from Black Forest Labs made with Pruna. We achieve a ~3x speedup over the original model with minimal quality loss.
Return CLIP features for the clip-vit-large-patch14 model
Practical face restoration algorithm for *old photos* or *AI-generated faces*
High resolution image Upscaler and Enhancer. Use at ClarityAI.co. A free Magnific alternative. Twitter/X: @philz1337x
Latest models
Source: Arc53/docsgpt-7b-mistral ✦ Quant: TheBloke/docsgpt-7B-mistral-AWQ ✦ DocsGPT is optimized for Documentation (RAG), fine-tuned for providing answers that are based on context
GFPGAN aims at developing Practical Algorithms for Real-world Face Restoration
Better than SDXL at both prompt adherence and image quality, by dataautogpt3
this is the replicate version of singing_voice_conversion from amphion
Animagine XL 2.0 is an advanced latent text-to-image diffusion model designed to create high-resolution, detailed anime images.
ElasticDiffusion: Training-free Arbitrary Size Image Generation
Super High Quality Depth Maps 🗺️: An End-to-End Tile-Based Framework 🏗️ for High-Resolution Monocular Metric Depth Estimation 🔍📏
A unique fusion that showcases exceptional prompt adherence and semantic understanding, it seems to be a step above base SDXL and a step closer to DALLE-3 in terms of prompt comprehension
Space-Time Diffusion Features for Zero-Shot Text-Driven Motion Transfer
Nous Hermes 2 - Yi-34B is a state of the art Yi Fine-tune, fine tuned on GPT-4 generated synthetic data
Terminus XL Otaku is a latent diffusion model that uses zero-terminal SNR noise schedule and velocity prediction objective at training and inference time.
works with inpainting and multi-controlnet + single-controlnet || ip-adapter + without ip adapter
Terminus XL Gamma is a new state-of-the-art latent diffusion model that uses zero-terminal SNR noise schedule and velocity prediction objective at training and inference time.
Source: SuperAGI/SAM ✦ Quant: TheBloke/SAM-AWQ ✦ SAM (Small Agentic Model), a 7B model that demonstrates impressive reasoning abilities despite its smaller size
Multi-controlnet, lora loading, img2img, inpainting
Try out akx/Poro-34B-gguf, Q5_K, This is 1000B checkpoint model
Amphion Singing Voice Conversion: DiffWaveNetSVC
Amazing photorealism with RealVisXL_V3.0, based on SDXL, trainable
Cog implementation of mir-aidj(Taejun Kim)'s 'All-In-One Music Structure Analyzer'
(Research only) IP-Adapter-FaceID can generate various style images conditioned on a face with only text prompts
DreamShaper is a general purpose SD model that aims at doing everything well, photos, art, anime, manga. It's designed to go against other general purpose models and pipelines like Midjourney and DALL-E.
DPO-SDXL Canny controlnet with LoRA support.
DreamShaper is a general purpose SD model that aims at doing everything well, photos, art, anime, manga. It's designed to match Midjourney and DALL-E.
Direct Preference Optimization (DPO) is a method to align diffusion models to text human preferences by directly optimizing on human comparison data
FacebookResearch/SeamlessM4T v2 - Massively Multilingual & Multimodal Machine Translation
The "Overall Best Performing Open Source 7B Model" for Coding + Generalization or Mathematical Reasoning
Source: kaist-ai/prometheus-13b-v1.0 ✦ Quant: TheBloke/prometheus-13B-v1.0-AWQ ✦ An alternative to GPT-4 when evaluating LLMs & Reward models for RLHF
Source: OpenBuddy/openbuddy-zephyr-7b-v14.1 ✦ Quant: TheBloke/openbuddy-zephyr-7B-v14.1-AWQ ✦ Open Multilingual Chatbot
AnimateDiff v3 + SparseCtrl: Animate Your Personalized Text-to-Image Diffusion Models without Specific Tuning. Created with Shimmer.
Blue Pencil XL v2 Model (Text2Img, Img2Img and Inpainting)
SDXL image generation using ComfyUI with LoRA trained on DreamBooth method.
Source: upstage/SOLAR-10.7B-Instruct-v1.0 ✦ Quant: TheBloke/SOLAR-10.7B-Instruct-v1.0-AWQ ✦ Elevating Performance with Upstage Depth UP Scaling!
AI-driven audio enhancement for your audio files, powered by Resemble AI
Zero-shot speech synthesizer for text-to-speech and voice conversion